# Anime character generation

A text-to-image model based on Stable Diffusion XL, specializing in anime style with exceptional handling of character details such as hands and feet, and delivering extreme lighting effects.

A low-rank adaptation generation model based on the Flux foundation model, specializing in generating semi-realistic style images

An anime-style image generation model enhanced based on the Animagine XL 3.1 model, optimized for high-quality anime character art generation.

LoRA model trained on 150 GPT-4 visually annotated WLOP artwork images for generating WLOP-style images

Animagine XL 3.1 is an upgraded version of the Animagine XL V3 series. It has been optimized on the basis of the previous work and can generate higher-quality anime images.

This is a test model highly similar to DALL·E 3, designed for text-to-image generation tasks.

MagicalMix_v1 is a text-to-image model based on Stable Diffusion, specializing in generating illustrations with a matte texture.

MoeDiffusion++ is a text-to-image model specialized in generating anime characters with black hair and twin tails, optimized for facial expressions and style control through the fusion of multiple base models.

Anything3.0 is an overfitted model, excelling in generating character images and specific details, capable of producing high-quality images even with suboptimal prompts.

This is a DreamBooth model trained on Waifu-Diffusion1.2, specifically designed to generate images resembling the character Aikawa Kina from 'Magia Record'.